github huggingface/diffusers v0.7.0
v0.7.0: Optimized for Apple Silicon, Improved Performance, Awesome Community
on GitHub

latest releases: v0.30.3, v0.30.2, v0.30.1...
22 months ago

❤️ PyTorch + Accelerate

⚠️ The PyTorch pipelines now require accelerate for improved model loading times!
Install Diffusers with pip install --upgrade diffusers[torch] to get everything in a single command.

🍎 Apple Silicon support with PyTorch 1.13

PyTorch and Apple have been working on improving mps support in PyTorch 1.13, so Apple Silicon is now a first-class citizen in diffusers 0.7.0!

Requirements

  • Mac computer with Apple silicon (M1/M2) hardware.
  • macOS 12.6 or later (13.0 or later recommended, as support is even better).
  • arm64 version of Python.
  • PyTorch 1.13.0 official release, installed from pip or the conda channels.

Memory efficient generation

Memory management is crucial to achieve fast generation speed. We recommend to always use attention slicing on Apple Silicon, as it drastically reduces memory pressure and prevents paging or swapping. This is especially important for computers with less than 64 GB of Unified RAM, and may be the difference between generating an image in seconds rather than in minutes. Use it like this: 

from diffusers import StableDiffusionPipeline

pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
pipe = pipe.to("mps")

# Recommended if your computer has < 64 GB of RAM
pipe.enable_attention_slicing()

prompt = "a photo of an astronaut riding a horse on mars"

# First-time "warmup" pass
_ = pipe(prompt, num_inference_steps=1)

image = pipe(prompt).images[0]
image.save("astronaut.png")

Continuous Integration

Our automated tests now include a full battery of tests on the mps device. This will be helpful to identify issues early and ensure the quality on Apple Silicon going forward.

See more details in the documentation.

💃 Dance Diffusion

diffusers goes audio 🎵 Dance Diffusion by Harmonai is the first audio model in 🧨Diffusers!

Try it out to generate some random music: 

from diffusers import DiffusionPipeline
import scipy

model_id = "harmonai/maestro-150k"
pipeline = DiffusionPipeline.from_pretrained(model_id)
pipeline = pipeline.to("cuda")

audio = pipeline(audio_length_in_s=4.0).audios[0]

# To save locally
scipy.io.wavfile.write("maestro_test.wav", pipe.unet.sample_rate, audio.transpose())

🎉 Euler schedulers

These are the Euler schedulers, from the paper Elucidating the Design Space of Diffusion-Based Generative Models by Karras et al. (2022). The diffusers implementation is based on the original k-diffusion implementation by Katherine Crowson. The Euler schedulers are fast, often times generating really good outputs with 20-30 steps. 

from diffusers import StableDiffusionPipeline, EulerDiscreteScheduler

euler_scheduler = EulerDiscreteScheduler.from_config("runwayml/stable-diffusion-v1-5", subfolder="scheduler")
pipeline = StableDiffusionPipeline.from_pretrained(
    "runwayml/stable-diffusion-v1-5", scheduler=euler_scheduler, revision="fp16", torch_dtype=torch.float16
)
pipeline.to("cuda")

prompt = "a photo of an astronaut riding a horse on mars"
image = pipeline(prompt, num_inference_steps=25).images[0]
from diffusers import StableDiffusionPipeline, EulerAncestralDiscreteScheduler

euler_ancestral_scheduler = EulerAncestralDiscreteScheduler.from_config("runwayml/stable-diffusion-v1-5", subfolder="scheduler")
pipeline = StableDiffusionPipeline.from_pretrained(
    "runwayml/stable-diffusion-v1-5", scheduler=euler_scheduler, revision="fp16", torch_dtype=torch.float16
)
pipeline.to("cuda")

prompt = "a photo of an astronaut riding a horse on mars"
image = pipeline(prompt, num_inference_steps=25).images[0]

🔥 Up to 2x faster inference with memory_efficient_attention

Even faster and memory efficient stable diffusion using the efficient flash attention implementation from xformers

To leverage it just make sure you have:

  • PyTorch > 1.12
  • Cuda available
  • Installed the xformers library 
from diffusers import StableDiffusionPipeline
import torch

pipe = StableDiffusionPipeline.from_pretrained(
    "runwayml/stable-diffusion-v1-5",
    revision="fp16",
    torch_dtype=torch.float16,
).to("cuda")

pipe.enable_xformers_memory_efficient_attention()

with torch.inference_mode():
    sample = pipe("a small cat")

# optional: You can disable it via
# pipe.disable_xformers_memory_efficient_attention()

🚀 Much faster loading

Thanks to accelerate, pipeline loading is much, much faster. There are two parts to it:

  • First, when a model is created PyTorch initializes its weights by default. This takes a good amount of time. Using low_cpu_mem_usage (enabled by default), no initialization will be performed.
  • Optionally, you can also use device_map="auto" to automatically select the best device(s) where the pre-trained weights will be initially sent to.

In our tests, loading time was more than halved on CUDA devices, and went down from 12s to 4s on an Apple M1 computer.

As a side effect, CPU usage will be greatly reduced during loading, because no temporary copies of the weights are necessary.

This feature requires PyTorch 1.9 or better and accelerate 0.8.0 or higher.

🎨 RePaint

RePaint allows to reuse any pretrained DDPM model for free-form inpainting by adding restarts to the denoising schedule. Based on the paper RePaint: Inpainting using Denoising Diffusion Probabilistic Models by Andreas Lugmayr et al. 

from diffusers import RePaintPipeline, RePaintScheduler

# Load the RePaint scheduler and pipeline based on a pretrained DDPM model
scheduler = RePaintScheduler.from_config("google/ddpm-ema-celebahq-256")
pipe = RePaintPipeline.from_pretrained("google/ddpm-ema-celebahq-256", scheduler=scheduler)
pipe = pipe.to("cuda")

generator = torch.Generator(device="cuda").manual_seed(0)
output = pipe(
    original_image=original_image,
    mask_image=mask_image,
    num_inference_steps=250,
    eta=0.0,
    jump_length=10,
    jump_n_sample=10,
    generator=generator,
)
inpainted_image = output.images[0]

image

🌍 Community Pipelines

Long Prompt Weighting Stable Diffusion

The Pipeline lets you input prompt without 77 token length limit. And you can increase words weighting by using "()" or decrease words weighting by using "[]". The Pipeline also lets you use the main use cases of the stable diffusion pipeline in a single class.
For a code example, see Long Prompt Weighting Stable Diffusion

  • [Community Pipelines] Long Prompt Weighting Stable Diffusion Pipelines by @SkyTNT in #907

Speech to Image

Generate an image from an audio sample using pre-trained OpenAI whisper-small and Stable Diffusion.
For a code example, see Speech to Image

Wildcard Stable Diffusion

A minimal implementation that allows for users to add "wildcards", denoted by __wildcard__ to prompts that are used as placeholders for randomly sampled values given by either a dictionary or a .txt file.
For a code example, see Wildcard Stable Diffusion

Composable Stable Diffusion

Use logic operators to do compositional generation.
For a code example, see Composable Stable Diffusion

  • Add Composable diffusion to community pipeline examples by @MarkRich in #951

Imagic Stable Diffusion

Image editing with Stable Diffusion.
For a code example, see Imagic Stable Diffusion

Seed Resizing

Allows to generate a larger image while keeping the content of the original image.
For a code example, see Seed Resizing

📝 Changelog

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